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Controllable Generation: Pix2Pix0, Attend and Excite, SEGA, SAG, ...

20 Feb 09:42
@patrickvonplaten patrickvonplaten
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Controllable Generation: Pix2Pix0, Attend and Excite, SEGA, SAG, ...

🎯 Controlling Generation

There has been much recent work on fine-grained control of diffusion networks!

Diffusers now supports:

  1. Instruct Pix2Pix
  2. Pix2Pix 0, more details in docs
  3. Attend and excite, more details in docs
  4. Semantic guidance, more details in docs
  5. Self-attention guidance, more details in docs
  6. Depth2image
  7. MultiDiffusion panorama, more details in docs

See our doc on controlling image generation and the individual pipeline docs for more details on the individual methods.

🆙 Latent Upscaler

Latent Upscaler is a diffusion model that is designed explicitly for Stable Diffusion. You can take the generated latent from Stable Diffusion and pass it into the upscaler before decoding with your standard VAE. Or you can take any image, encode it into the latent space, use the upscaler, and decode it. It is incredibly flexible and can work with any SD checkpoints.

Original output image 2x upscaled output image

The model was developed by Katherine Crowson in collaboration with Stability AI

from diffusers import StableDiffusionLatentUpscalePipeline, StableDiffusionPipeline
import torch

pipeline = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16)
pipeline.to("cuda")

upscaler = StableDiffusionLatentUpscalePipeline.from_pretrained("stabilityai/sd-x2-latent-upscaler", torch_dtype=torch.float16)
upscaler.to("cuda")

prompt = "a photo of an astronaut high resolution, unreal engine, ultra realistic"
generator = torch.manual_seed(33)

# we stay in latent space! Let's make sure that Stable Diffusion returns the image
# in latent space
low_res_latents = pipeline(prompt, generator=generator, output_type="latent").images

upscaled_image = upscaler(
    prompt=prompt,
    image=low_res_latents,
    num_inference_steps=20,
    guidance_scale=0,
    generator=generator,
).images[0]

# Let's save the upscaled image under "upscaled_astronaut.png"
upscaled_image.save("astronaut_1024.png")

# as a comparison: Let's also save the low-res image
with torch.no_grad():
    image = pipeline.decode_latents(low_res_latents)
image = pipeline.numpy_to_pil(image)[0]

image.save("astronaut_512.png")

⚡ Optimization

In addition to new features and an increasing number of pipelines, diffusers cares a lot about performance. This release brings a number of optimizations that you can turn on easily.

xFormers

Memory efficient attention, as implemented by xFormers, has been available in diffusers for some time. The problem was that installing xFormers could be complicated because there were no official pip wheels (or they were outdated), and you had to resort to installing from source.

From xFormers 0.0.16, official pip wheels are now published with every release, so installing and using xFormers is now as simple as these two steps:

  1. pip install xformers in your terminal.
  2. pipe.enable_xformers_memory_efficient_attention() in your code to opt-in in your pipelines.

These actions will unlock dramatic memory savings, and usually faster inference too!

See more details in the documentation.

Torch 2.0

Speaking of memory-efficient attention, Accelerated PyTorch 2.0 Transformers now comes with built-in native support for it! When PyTorch 2.0 is released you'll no longer have to install xFormers or any third-party package to take advantage of it. In diffusers we are already preparing for that, and it works out of the box. So, if you happen to be using the latest "nightlies" of PyTorch 2.0 beta, then you're all set – diffusers will use Accelerated PyTorch 2.0 Transformers by default.

In our tests, the built-in PyTorch 2.0 implementation is usually as fast as xFormers', and sometimes even faster. Performance depends on the card you are using and whether you run your code in float16 or float32, so check our documentation for details.

Coarse-grained CPU offload

Community member @keturn, with whom we have enjoyed thoughtful software design conversations, called our attention to the fact that enabling sequential cpu offloading via enable_sequential_cpu_offload worked great to save a lot of memory, but made inference much slower.

This is because enable_sequential_cpu_offload() is optimized for memory, and it recursively works across all the submodules contained in a model, moving them to GPU when they are needed and back to CPU when another submodule needs to run. These cpu-to-gpu-to-cpu transfers happen hundreds of times during the stable diffusion denoising loops, because the UNet runs multiple times and it consists of several PyTorch modules.

This release of diffusers introduces a coarser enable_model_cpu_offload() pipeline API, which copies whole models (not modules) to GPU and makes sure they stay there until another model needs to run. The consequences are:

  • Less memory savings than enable_sequential_cpu_offload, but:
  • Almost as fast inference as when the pipeline is used without any type of offloading.

Pix2Pix Zero

Remember the CycleGAN days where one would turn a horse into a zebra in an image while keeping the rest of the content almost untouched? Well, that day has arrived but in the context of Diffusion. Pix2Pix Zero allows users to edit a particular image (be it real or generated), targeting a source concept (horse, for example) and replacing it with a target concept (zebra, for example).

Input image Edited image
original edited

Pix2Pix Zero was proposed in Zero-shot Image-to-Image Translation. The StableDiffusionPix2PixZeroPipeline allows you to

  1. Edit an image generated from an input prompt
  2. Provide an input image and edit it

For the latter, it uses the newly introduced DDIMInverseScheduler to first obtain the inverted noise from the input image and use that in the subsequent generation process.

Both of the use cases leverage the idea of "edit directions", used for steering the generation toward the target concept gradually from the source concept. To know more, we recommend checking out the official documentation.

Attend and excite

Attend-and-Excite: Attention-Based Semantic Guidance for Text-to-Image Diffusion Models. Attend-and-Excite, guides the generative model to modify the cross-attention values during the image synthesis process to generate images that more faithfully depict the input text prompt. It allows creating images that are more semantically faithful with respect to the input text prompts. Thanks to community contributor @evinpinar for leading the charge to add this pipeline!

Semantic guidance

Semantic Guidance for Diffusion Models was proposed in SEGA: Instructing Diffusion using Semantic Dimensions and provides strong semantic control over image generation. Small changes to the text prompt usually result in entirely different output images. However, with SEGA, a variety of changes to the image are enabled that can be controlled easily and intuitively and stay true to the original image composition. Thanks to the lead author of SEFA, Manuel (@manuelbrack), who added the pipeline in #2223.

Here is a simple demo:

import torch
from diffusers import SemanticStableDiffusionPipeline

pipe = SemanticStableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe = pipe.to("cuda")

out = pipe(
    prompt="a photo of the face of a woman",
    num_images_per_prompt=1,
    guidance_scale=7,
    editing_prompt=[
        "smiling, smile",  # Concepts to apply
        "glasses, wearing glasses",
        "curls, wavy hair, curly hair",
        "beard, full beard, mustache",
    ],
    reverse_editing_direction=[False, False, False, False],  # Direction of guidance i.e. increase all concepts
    edit_warmup_steps=[10, 10, 10, 10],  # Warmup period for each concept
...
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kashif, keturn, and 49 other contributors
Assets 2

v0.12.1: Patch Release to fix local files only

27 Jan 07:56
@patrickvonplaten patrickvonplaten
v0.12.1
374ec16
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v0.12.1: Patch Release to fix local files only

Make sure cached models can be loaded in offline mode.

Contributors

patrickvonplaten
Assets 2

Instruct-Pix2Pix, DiT, LoRA

25 Jan 16:10
@patrickvonplaten patrickvonplaten
v0.12.0
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Instruct-Pix2Pix, DiT, LoRA

🪄 Instruct-Pix2Pix

Instruct-Pix2Pix is a Stable Diffusion model fine-tuned for editing images from human instructions. Given an input image and a written instruction that tells the model what to do, the model follows these instructions to edit the image.

image

The model was released with the paper InstructPix2Pix: Learning to Follow Image Editing Instructions. More information about the model can be found in the paper.

pip install diffusers transformers safetensors accelerate

import PIL
import requests
import torch
from diffusers import StableDiffusionInstructPix2PixPipeline

model_id = "timbrooks/instruct-pix2pix"
pipe = StableDiffusionInstructPix2PixPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")

url = "https://huggingface.co/datasets/diffusers/diffusers-images-docs/resolve/main/mountain.png"
def download_image(url):
    image = PIL.Image.open(requests.get(url, stream=True).raw)
    image = PIL.ImageOps.exif_transpose(image)
    image = image.convert("RGB")
    return image
image = download_image(url)

prompt = "make the mountains snowy"
edit = pipe(prompt, image=image, num_inference_steps=20, image_guidance_scale=1.5, guidance_scale=7).images[0]
images[0].save("snowy_mountains.png")

🤖 DiT

Diffusion Transformers (DiTs) is a class conditional latent diffusion model which replaces the commonly used U-Net backbone with a transformer operating on latent patches. The pretrained model is trained on the ImageNet-1K dataset and is able to generate class conditional images of 256x256 or 512x512 pixels.

dit

The model was released with the paper Scalable Diffusion Models with Transformers.

import torch
from diffusers import DiTPipeline

model_id = "facebook/DiT-XL-2-256"
pipe = DiTPipeline.from_pretrained(model_id, torch_dtype=torch.float16).to("cuda")

# pick words that exist in ImageNet
words = ["white shark", "umbrella"]
class_ids = pipe.get_label_ids(words)

output = pipe(class_labels=class_ids)
image = output.images[0]  # label 'white shark'

⚡ LoRA

LoRA is a technique for performing parameter-efficient fine-tuning for large models. LoRA works by adding so-called "update matrices" to specific blocks of a pre-trained model. During fine-tuning, only these update matrices are updated while the pre-trained model parameters are kept frozen. This allows us to achieve greater memory efficiency as well as easier portability during fine-tuning.

LoRA was proposed in LoRA: Low-Rank Adaptation of Large Language Models. In the original paper, the authors investigated LoRA for fine-tuning large language models like GPT-3. cloneofsimo was the first to try out LoRA training for Stable Diffusion in the popular lora GitHub repository.

Diffusers now supports LoRA! This means you can now fine-tune a model like Stable Diffusion using consumer GPUs like Tesla T4 or RTX 2080 Ti. LoRA support was added to UNet2DConditionModel and DreamBooth training script by @patrickvonplaten in #1884.

By using LoRA, the fine-tuned checkpoints will be just 3 MBs in size. After fine-tuning, you can use the LoRA checkpoints like so:

from diffusers import StableDiffusionPipeline
import torch

model_path = "sayakpaul/sd-model-finetuned-lora-t4"
pipe = StableDiffusionPipeline.from_pretrained("CompVis/stable-diffusion-v1-4", torch_dtype=torch.float16)
pipe.unet.load_attn_procs(model_path)
pipe.to("cuda")

prompt = "A pokemon with blue eyes."
image = pipe(prompt, num_inference_steps=30, guidance_scale=7.5).images[0]
image.save("pokemon.png")

pokemon-image

You can follow these resources to know more about how to use LoRA in diffusers:

📐 Customizable Cross Attention

LoRA leverages a new method to customize the cross attention layers deep in the UNet. This can be useful for other creative approaches such as Prompt-to-Prompt, and it makes it easier to apply optimizers like xFormers. This new "attention processor" abstraction was created by @patrickvonplaten in #1639 after discussing the design with the community, and we have used it to rewrite our xFormers and attention slicing implementations!

🌿 Flax => PyTorch

A long requested feature, prolific community member @camenduru took up the gauntlet in #1900 and created a way to convert Flax model weights for PyTorch. This means that you can train or fine-tune models super fast using Google TPUs, and then convert the weights to PyTorch for everybody to use. Thanks @camenduru!

🌀 Flax Img2Img

Another community member, @dhruvrnaik, ported the image-to-image pipeline to Flax in #1355! Using a TPU v2-8 (available in Colab's free tier), you can generate 8 images at once in a few seconds!

🎲 DEIS Scheduler

DEIS (Diffusion Exponential Integrator Sampler) is a new fast mult step scheduler that can generate high-quality samples in fewer steps.
The scheduler was introduced in the paper Fast Sampling of Diffusion Models with Exponential Integrator. More information about the scheduler can be found in the paper.

from diffusers import StableDiffusionPipeline, DEISMultistepScheduler
import torch

pipe = StableDiffusionPipeline.from_pretrained("runwayml/stable-diffusion-v1-5", torch_dtype=torch.float16)
pipe.scheduler = DEISMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")

prompt = "a photo of an astronaut riding a horse on mars"
generator = torch.Generator(device="cuda").manual_seed(0)
image = pipe(prompt, generator=generator, num_inference_steps=25).images[0
  • feat : add log-rho deis multistep scheduler by @qsh-zh #1432

Reproducibility

One can now pass CPU generators to all pipelines even if the pipeline is on GPU. This ensures
much better reproducibility across GPU hardware:

import torch
from diffusers import DDIMPipeline
import numpy as np

model_id = "google/ddpm-cifar10-32"

# load model and scheduler
ddim = DDIMPipeline.from_pretrained(model_id)
ddim.to("cuda")

# create a generator for reproducibility
generator = torch.manual_seed(0)

# run pipeline for just two steps and return numpy tensor
image = ddim(num_inference_steps=2, output_type="np", generator=generator).images
print(np.abs(image).sum())

See: #1902 and https://huggingface.co/docs/diffusers/using-diffusers/reproducibility

Important New Guides

Important Bug Fixes

  • Don't download safetensors if library is not installed: #2057
  • Make sure that save_pretrained(...) doesn't accidentally delete files: #2038
  • Fix CPU offload docs for maximum memory gain: #1968
  • Fix conversion for exotically sorted weight names: #1959
  • Fix intermediate checkpointing for textual inversion, thanks @lstein #2072

All commits

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kig, kashif, and 53 other contributors
Assets 2

v0.11.1: Patch release

20 Dec 16:32
@anton-l anton-l
v0.11.1
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v0.11.1: Patch release

This patch release fixes a bug with num_images_per_prompt in the UnCLIPPipeline

Contributors

patil-suraj
Assets 2

v0.11.0: Karlo UnCLIP, safetensors, pipeline versions

19 Dec 17:47
@anton-l anton-l
v0.11.0
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v0.11.0: Karlo UnCLIP, safetensors, pipeline versions

🪄 Karlo UnCLIP by Kakao Brain

Karlo is a text-conditional image generation model based on OpenAI's unCLIP architecture with the improvement over the standard super-resolution model from 64px to 256px, recovering high-frequency details in a small number of denoising steps.

This alpha version of Karlo is trained on 115M image-text pairs, including COYO-100M high-quality subset, CC3M, and CC12M.
For more information about the architecture, see the Karlo repository: https://github.com/kakaobrain/karlo
image

pip install diffusers transformers safetensors accelerate

import torch
from diffusers import UnCLIPPipeline

pipe = UnCLIPPipeline.from_pretrained("kakaobrain/karlo-v1-alpha", torch_dtype=torch.float16)
pipe = pipe.to("cuda")

prompt = "a high-resolution photograph of a big red frog on a green leaf."
image = pipe(prompt).images[0]

img

:octocat: Community pipeline versioning

The community pipelines hosted in diffusers/examples/community will now follow the installed version of the library.

E.g. if you have diffusers==0.9.0 installed, the pipelines from the v0.9.0 branch will be used: https://github.com/huggingface/diffusers/tree/v0.9.0/examples/community

If you've installed diffusers from source, e.g. with pip install git+https://github.com/huggingface/diffusers then the latest versions of the pipelines will be fetched from the main branch.

To change the custom pipeline version, set the custom_revision variable like so:

pipeline = DiffusionPipeline.from_pretrained(
     "google/ddpm-cifar10-32", custom_pipeline="one_step_unet", custom_revision="0.10.2"
)

🦺 safetensors

Many of the most important checkpoints now have https://github.com/huggingface/safetensors available. Upon installing safetensors with:

pip install safetensors

You will see a nice speed-up when loading your model 🚀

Some of the most improtant checkpoints have safetensor weights added now:

Batched generation bug fixes 🐛

We fixed a lot of bugs for batched generation. All pipelines should now correctly process batches of prompts and images 🤗
Also we made it much easier to tweak images with reproducible seeds:
https://huggingface.co/docs/diffusers/using-diffusers/reusing_seeds

📝 Changelog

Contributors

maruel, msiedlarek, and 22 other contributors
Assets 2

v0.10.2: Patch release

09 Dec 18:19
@anton-l anton-l
v0.10.2
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v0.10.2: Patch release

This patch removes the hard requirement for transformers>=4.25.1 in case external libraries were downgrading the library upon startup in a non-controllable way.

🚨🚨🚨 Note that xformers in not automatically enabled anymore 🚨🚨🚨

The reasons for this are given here: #1640 (comment):

We should not automatically enable xformers for three reasons:

It's not PyTorch-like API. PyTorch doesn't by default enable all the fastest options available
We allocate GPU memory before the user even does .to("cuda")
This behavior is not consistent with cases where xformers is not installed

=> This means: If you were used to have xformers automatically enabled, please make sure to add the following now:

from diffusers.utils.import_utils import is_xformers_available

unet = ... # load unet

if is_xformers_available():
    try:
        unet.enable_xformers_memory_efficient_attention(True)
    except Exception as e:
        logger.warning(
            "Could not enable memory efficient attention. Make sure xformers is installed"
            f" correctly and a GPU is available: {e}"
        )

for the UNet (e.g. in dreambooth) or for the pipeline:

from diffusers.utils.import_utils import is_xformers_available

pipe = ... # load pipeline

if is_xformers_available():
    try:
        pipe.enable_xformers_memory_efficient_attention(True)
    except Exception as e:
        logger.warning(
            "Could not enable memory efficient attention. Make sure xformers is installed"
            f" correctly and a GPU is available: {e}"
        )

Contributors

patrickvonplaten and anton-l
Assets 2

v0.10.1: Patch release

09 Dec 14:12
@anton-l anton-l
v0.10.1
5bc6cd5
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v0.10.1: Patch release

This patch returns enable_xformers_memory_efficient_attention() to UNet2DCondition to restore backward compatibility.

Contributors

patrickvonplaten
Assets 2

v0.10.0: Depth Guidance and Safer Checkpoints

08 Dec 18:32
@anton-l anton-l
v0.10.0
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v0.10.0: Depth Guidance and Safer Checkpoints

🐳 Depth-Guided Stable Diffusion and 2.1 checkpoints

The new depth-guided stable diffusion model is fully supported in this release. The model is conditioned on monocular depth estimates inferred via MiDaS and can be used for structure-preserving img2img and shape-conditional synthesis.

image

Installing the transformers library from source is required for the MiDaS model:

pip install --upgrade git+https://github.com/huggingface/transformers/
import torch
import requests
from PIL import Image
from diffusers import StableDiffusionDepth2ImgPipeline

pipe = StableDiffusionDepth2ImgPipeline.from_pretrained(
   "stabilityai/stable-diffusion-2-depth",
   torch_dtype=torch.float16,
).to("cuda")

url = "http://images.cocodataset.org/val2017/000000039769.jpg"
init_image = Image.open(requests.get(url, stream=True).raw)

prompt = "two tigers"
n_propmt = "bad, deformed, ugly, bad anotomy"
image = pipe(prompt=prompt, image=init_image, negative_prompt=n_propmt, strength=0.7).images[0]

The updated Stable Diffusion 2.1 checkpoints are also released and fully supported:

🦺 Safe Tensors

We now support SafeTensors: a new simple format for storing tensors safely (as opposed to pickle) that is still fast (zero-copy).

Format Safe Zero-copy Lazy loading No file size limit Layout control Flexibility Bfloat16
pickle (PyTorch) ✗ ✗ ✗ ✓ ✗ ✓ ✓
H5 (Tensorflow) ✓ ✗ ✓ ✓ ~ ~ ✗
SavedModel (Tensorflow) ✓ ✗ ✗ ✓ ✓ ✗ ✓
MsgPack (flax) ✓ ✓ ✗ ✓ ✗ ✗ ✓
SafeTensors ✓ ✓ ✓ ✓ ✓ ✗ ✓

**More details about the comparison here: https://github.com/huggingface/safetensors#yet-another-format-

pip install safetensors

from diffusers import StableDiffusionPipeline

pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1")
pipe.save_pretrained("./safe-stable-diffusion-2-1", safe_serialization=True)

# you can also push this checkpoint to the HF Hub and load from there
safe_pipe = StableDiffusionPipeline.from_pretrained("./safe-stable-diffusion-2-1")

New Pipelines

🖌️ Paint-by-example

An implementation of Paint by Example: Exemplar-based Image Editing with Diffusion Models by Binxin Yang, Shuyang Gu, Bo Zhang, Ting Zhang, Xuejin Chen, Xiaoyan Sun, Dong Chen, Fang Wen

image

import PIL
import requests
import torch
from io import BytesIO
from diffusers import DiffusionPipeline

def download_image(url):
    response = requests.get(url)
    return PIL.Image.open(BytesIO(response.content)).convert("RGB")

img_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/image/example_1.png"
mask_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/mask/example_1.png"
example_url = "https://raw.githubusercontent.com/Fantasy-Studio/Paint-by-Example/main/examples/reference/example_1.jpg"

init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))
example_image = download_image(example_url).resize((512, 512))

pipe = DiffusionPipeline.from_pretrained("Fantasy-Studio/Paint-by-Example", torch_dtype=torch.float16)
pipe = pipe.to("cuda")

image = pipe(image=init_image, mask_image=mask_image, example_image=example_image).images[0]

Audio Diffusion and Latent Audio Diffusion

Audio Diffusion leverages the recent advances in image generation using diffusion models by converting audio samples to and from mel spectrogram images.

from IPython.display import Audio
from diffusers import DiffusionPipeline

pipe = DiffusionPipeline.from_pretrained("teticio/audio-diffusion-ddim-256").to("cuda")

output = pipe()
display(output.images[0])
display(Audio(output.audios[0], rate=pipe.mel.get_sample_rate()))

[Experimental] K-Diffusion pipeline for Stable Diffusion

This pipeline is added to support the latest schedulers from @crowsonkb's k-diffusion
The purpose of this pipeline is to compare scheduler implementations and updates, so new features from other pipelines are unlikely to be supported!

pip install k-diffusion

from diffusers import StableDiffusionKDiffusionPipeline
import torch

pipe = StableDiffusionKDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1-base")
pipe = pipe.to("cuda")

pipe.set_scheduler("sample_heun")
image = pipe("astronaut riding horse", num_inference_steps=25).images[0]

New Schedulers

Heun scheduler inspired by Karras et. al

Algorithm 1 of Karras et. al. Scheduler ported from @crowsonkb’s k-diffusion

from diffusers import HeunDiscreteScheduler

pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1")
pipe.scheduler = HeunDiscreteScheduler.from_config(pipe.scheduler.config)

Single step DPM-Solver

Original paper can be found here and the improved version. The original implementation can be found here.

  • Add Singlestep DPM-Solver (singlestep high-order schedulers) by @LuChengTHU in #1442
from diffusers import DPMSolverSinglestepScheduler

pipe = StableDiffusionPipeline.from_pretrained("stabilityai/stable-diffusion-2-1")
pipe.scheduler = DPMSolverSinglestepScheduler.from_config(pipe.scheduler.config)

📝 Changelog

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Contributors

kig, MatthieuBizien, and 27 other contributors
Assets 2

v0.9.0: Stable Diffusion 2

25 Nov 17:17
@anton-l anton-l
v0.9.0
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v0.9.0: Stable Diffusion 2

🎨 Stable Diffusion 2 is here!

Installation

pip install diffusers[torch]==0.9 transformers

Stable Diffusion 2.0 is available in several flavors:

Stable Diffusion 2.0-V at 768x768

New stable diffusion model (Stable Diffusion 2.0-v) at 768x768 resolution. Same number of parameters in the U-Net as 1.5, but uses OpenCLIP-ViT/H as the text encoder and is trained from scratch. SD 2.0-v is a so-called v-prediction model.

image

import torch
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler

repo_id = "stabilityai/stable-diffusion-2"
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")

prompt = "High quality photo of an astronaut riding a horse in space"
image = pipe(prompt, guidance_scale=9, num_inference_steps=25).images[0]
image.save("astronaut.png")

Stable Diffusion 2.0-base at 512x512

The above model is finetuned from SD 2.0-base, which was trained as a standard noise-prediction model on 512x512 images and is also made available.

image

import torch
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler

repo_id = "stabilityai/stable-diffusion-2-base"
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")

prompt = "High quality photo of an astronaut riding a horse in space"
image = pipe(prompt, num_inference_steps=25).images[0]
image.save("astronaut.png")

Stable Diffusion 2.0 for Inpanting

This model for text-guided inpanting is finetuned from SD 2.0-base. Follows the mask-generation strategy presented in LAMA which, in combination with the latent VAE representations of the masked image, are used as an additional conditioning.

image

import PIL
import requests
import torch
from io import BytesIO
from diffusers import DiffusionPipeline, DPMSolverMultistepScheduler

def download_image(url):
    response = requests.get(url)
    return PIL.Image.open(BytesIO(response.content)).convert("RGB")

img_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo.png"
mask_url = "https://raw.githubusercontent.com/CompVis/latent-diffusion/main/data/inpainting_examples/overture-creations-5sI6fQgYIuo_mask.png"
init_image = download_image(img_url).resize((512, 512))
mask_image = download_image(mask_url).resize((512, 512))

repo_id = "stabilityai/stable-diffusion-2-inpainting"
pipe = DiffusionPipeline.from_pretrained(repo_id, torch_dtype=torch.float16, revision="fp16")
pipe.scheduler = DPMSolverMultistepScheduler.from_config(pipe.scheduler.config)
pipe = pipe.to("cuda")

prompt = "Face of a yellow cat, high resolution, sitting on a park bench"
image = pipe(prompt=prompt, image=init_image, mask_image=mask_image, num_inference_steps=25).images[0]
image.save("yellow_cat.png")

Stable Diffusion X4 Upscaler

The model was trained on crops of size 512x512 and is a text-guided latent upscaling diffusion model. In addition to the textual input, it receives a noise_level as an input parameter, which can be used to add noise to the low-resolution input according to a predefined diffusion schedule.

image

import requests
from PIL import Image
from io import BytesIO
from diffusers import StableDiffusionUpscalePipeline
import torch

model_id = "stabilityai/stable-diffusion-x4-upscaler"
pipeline = StableDiffusionUpscalePipeline.from_pretrained(model_id, revision="fp16", torch_dtype=torch.float16)
pipeline = pipeline.to("cuda")

url = "https://huggingface.co/datasets/hf-internal-testing/diffusers-images/resolve/main/sd2-upscale/low_res_cat.png"
response = requests.get(url)
low_res_img = Image.open(BytesIO(response.content)).convert("RGB")
low_res_img = low_res_img.resize((128, 128))

prompt = "a white cat"
upscaled_image = pipeline(prompt=prompt, image=low_res_img).images[0]
upscaled_image.save("upsampled_cat.png")

Saving & Loading is fixed for Versatile Diffusion

Previously there was a 🐛 when saving & loading versatile diffusion - this is fixed now so that memory efficient saving & loading works as expected.

📝 Changelog

Contributors

kashif, pcuenca, and 3 other contributors
Assets 2

v0.8.1: Patch release

24 Nov 00:09
@anton-l anton-l
v0.8.1
3876b20
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v0.8.1: Patch release

This patch release fixes an error with CLIPVisionModelWithProjection imports on a non-git transformers installation.

⚠️ Please upgrade with pip install --upgrade diffusers or pip install diffusers==0.8.1

Contributors

patrickvonplaten
Assets 2